# Text-to-Image Generation

A ControlNet PEFT LoRA model based on terminusresearch/auraflow-v0.3 for high-quality text-to-image conversion

MMaDA is a novel class of multimodal diffusion foundation models, excelling in various domains such as text reasoning, multimodal understanding, and text-to-image generation.

A diffusion-based text-to-image generation model capable of producing realistic images from text inputs

A stable diffusion model implemented with JAX and Flax, focusing on text-to-image generation tasks.

Illustrious-XL-v2.0 is a foundational model for text-to-image generation, suitable for various image generation tasks.

HiDream-I1-Full is a GGUF-format text-to-image generation model designed for image generation tasks.

HiDream-I1-Dev is an image generation model based on GGUF format conversion, supporting text-to-image generation tasks.

A text-to-image generation model based on black-forest-labs/FLUX.1-dev foundation model, fine-tuned using Low-Rank Adaptation (LoRA) technology

This is a LoRA fine-tuned model based on FLUX.1-dev, specifically designed for generating images in children's book style.

LoRA adapter trained on FLUX.1-dev model for text-to-image generation

LoRA fine-tuned model based on FLUX.1-dev, specializing in text-to-image generation tasks

XPEN is a text-to-image generation model based on diffusion models and LoRA technology

The Flux model employs a 4-bit Transformer and T5 encoder for text-to-image generation tasks, supporting non-commercial use.

LoRA fine-tuned model based on FLUX.1-dev, specializing in text-to-image generation tasks with Midjourney style integration

A Face ID customization model trained on the FaceID-6M dataset, capable of generating customized facial images based on text descriptions.

Janus-Pro is a novel autoregressive framework that unifies multimodal understanding and generation capabilities. By decoupling visual encoding paths, it uses a single Transformer architecture to handle multimodal tasks.

A text-to-image generation model based on Black Forest Lab's FLUX.1 development version, specializing in generating Pepe the Frog-related image content.

Lumina-Image 2.0 is a flow-based diffusion transformer model with 2 billion parameters, specializing in text-to-image generation tasks.

The Flex.1 alpha 8B model is a text-to-image generation model converted to GGUF format, specifically designed for the ComfyUI-GGUF tool.

A GLIGEN adapter supporting the SDXL version, providing HuggingFace-style pipelines for text-to-image generation tasks.

This is a direct GGUF conversion of the flux.1-lite-8B model, designed for the ComfyUI-GGUF framework.

FLUX.1 [Development Version] is a 12-billion-parameter rectified flow transformer model capable of generating high-quality images from text descriptions, featuring exceptional prompt-following capabilities and an open-weight design.

This fine-tuned model optimizes short prompts into detailed ones, enabling image generation based on Stable Diffusion or Flux technology.

A text-to-image generation model based on fal/AuraFlow-v0.3, offering GGUF quantized and FP8 scaled versions, suitable for ComfyUI environments.

A text-to-image generation model based on LoRA and diffusion models, capable of producing high-quality images from text prompts.

LoRA adapter trained on FLUX.1-dev model for text-to-image generation

Stable Diffusion is a text-to-image diffusion model capable of generating high-quality images from text descriptions.

This is a LoRA adaptation weight for stabilityai/stable-diffusion-xl-base-1.0, specifically designed for generating images of LEO cats.

A LoRA model based on FLUX.1-dev, specifically designed for generating text images with condensation effects on windows

Flux-Mini is a 3.2B-parameter MMDiT model distilled from Flux-dev, focusing on efficient text-to-image generation tasks.

FLUX.1-Fill-dev is a text-to-image generation model, specializing in image inpainting tasks.

FLUX.1-Redux-dev is a text-to-image generation model focused on converting textual descriptions into high-quality images.

FLUX.1-Redux-dev is a text-to-image generation model based on the FLUX technology stack, supporting the English language and adopting a non-commercial license.

A text-to-image generation model based on Stable Diffusion XL (SDXL), supporting ControlNet functionality.

FLUX.1-Canny-dev is a text-to-image generation model focused on image generation tasks.

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from text inputs.

A text-to-image generation model based on Stable Diffusion XL, supporting ControlNet functionality.

A text-to-image generation model based on Stable Diffusion XL, supporting ControlNet functionality

GGUF quantized version of Stable Diffusion 3.5 Medium, supporting text-to-image generation tasks.

A text-to-image generation model based on Stable Diffusion XL, supporting fine-grained image generation through ControlNet.